Amazing that their examples at the bottom of the page still show really messed up human hands.
I've personally observed that the drawing of hands in Midjourney and SD has been getting incrementally better release after release.
Stable Diffusion XL 1.0
91–100 of 182 posts
Re: Stable Diffusion XL 1.0
#92Earlier quoted context omitted.
You get access to the weights instantly if you apply for them. It's basically not a hurdle. (I've been having fun with this for a few days. https://huggingface.co/stabilityai/stable-diffusion-xl-base-... Not sure there's much of a difference with the 1.0 version.)
The ones you can apply for access to are the 0.9 weights, which have been available for a couple of weeks. Unless the SDXL 1.0 weights are also available by application somewhere that I'm unaware.
Re: Stable Diffusion XL 1.0
#93I'm out of date on the image-generating side of AI, but I'd like to check things out. What's the best tool for image generation that's available on a website right now? Ie, not a model that I have to run locally.
If you want to play around with Stable Diffusion XL: https://clipdrop.co
I just tried this and the UI is very nice (better than dreamstudio), with nice tool integration, and image quality is definitely going up with each new release. You can see a few results at fb.com/onlyrolydog (along with a lot of other canine nonsense).
Re: Stable Diffusion XL 1.0
#94Amazing that their examples at the bottom of the page still show really messed up human hands.
I've personally observed that the drawing of hands in Midjourney and SD has been getting incrementally better release after release.
That's why I'm amazed they picked images with totally borked up hands to put on their press release. Truth in advertising!
Re: Stable Diffusion XL 1.0
#95I hope someday there’s a version of this or something comparable to it that can run on <8gb consumer hardware. The main selling point of Stable Diffusion was its ability to run in that environment.
Give InvokeAI a try.
https://github.com/invoke-ai/InvokeAI
Edit: Spec required from the documentation
You will need one of the following:
An NVIDIA-based graphics card with 4 GB or more VRAM memory. 6-8 GB of VRAM is highly recommended for rendering using the Stable Diffusion XL models
An Apple computer with an M1 chip.
An AMD-based graphics card with 4GB or more VRAM memory (Linux only), 6-8 GB for XL rendering.Re: Stable Diffusion XL 1.0
#96Midjourney is still going to be hard to beat imo. Comparing SD to MJ is a little unfair considering their applications and flexibility, but I do really enjoy the "out of the box" experience that comes with MJ.
Re: Stable Diffusion XL 1.0
#97Earlier quoted context omitted.
https://playgroundai.com/create Not affiliated in anyway and not very involved in the space. I just wanted to generate some images a few weeks ago and was looking for somewhere I could do that for free. The link above lets you do that but I suggest you look up prompts because its a lot more involved than I expected.
Any particularly useful resources for looking into prompts?
I used this:
https://learnwithnaseem.com/best-playground-ai-prompts-for-a...
I just took the ones I liked and then deleted out the words that were specific to that image and left the ones that were providing the style of the image. So for example on the first one I would delete "an cute kitsune in florest" but would keep "colorfully fantast concept art". Then I just added a comma separated list of the of the features I wanted in my picture. It took a lot more trial and error than I thought and adding sentences seemed to be worse than just individual words. I am sure I barely scratched the surface of interfacing with the tool correctly but the space is moving so fast its not the kind of thing I want to spend my time learning right now just to have that knowledge deprecate in 6 months.
Re: Stable Diffusion XL 1.0
#98I always wondered why the vision models don't seem to be following the whole "scale up as much as possible" mantra that has defined the language models of the past few years (to the same extent). Even 3.5 billion parameters is absolutely nothing compared to the likes of GPT-3, 3.5, 4, or even the larger open-source language models (e.g. LLaMA-65B). Is it just an engineering challenge that no one has stepped up for ye…
Diffusion is more parameter-efficient and you quickly saturate the target fidelity, especially with some refiner cascade. It's a solved problem. You do not need more than maybe 4B total. Images are far more redundant than text. In fact, most interesting papers since Imagen show that you get more mileage out of scaling the text encoder part, which is, of course, a Transformer. This is what drives accuracy, text render…
> Images are far more redundant than text.
"A picture is worth a thousand words" - I wonder how (in)accurate this popular saying turned out to be? :D
Re: Stable Diffusion XL 1.0
#99Is there anything like this for the vector landscape?
This may just be due to the iterative denoising approach a lot of these models take but they only seem to work well when creating raster style images.
In my experience when you ask them to create logos, shirt designs, illustrations, they tend to not work as well and introduce a lot of artifacts, distortions, incorrect spellings etc.
Re: Stable Diffusion XL 1.0
#100Midjourney is still going to be hard to beat imo. Comparing SD to MJ is a little unfair considering their applications and flexibility, but I do really enjoy the "out of the box" experience that comes with MJ.
Different use case.
I can run SDXL 1.0 offline from my home. I can’t do this with Midjourney.
A closed source model that doesn’t have the limitation of running on consumer level GPUs will have certain advantages.